Although diffusion prior is rising as a powerful solution for blind face restoration (BFR), the inherent gap between the vanilla diffusion model and BFR settings hinders its seamless adaptation. The gap mainly stems from the discrepancy between 1) high-quality (HQ) and low-quality (LQ) images and 2) synthesized and real-world images.The vanilla diffusion model is trained on images with no or less degredations, while BFR handles moderately to severely degraded images.Additionally, LQ images used for training are synthesized by a naive degradation model with limited degradation patterns, which fails to simulate the complex and unknown degradations in real-world scenarios.In this work, we use a unified network FLIPNET that switches between two modes to address specific gaps.In restoration mode, the model gradually integrates BFR-oriented features and face embeddings from LQ images to achieve authentic and faithful face restoration.In degradation mode, the model synthesizes real-world like degraded images based on the knowledge learned from real-world degradation datasets.Extensive evaluations on benchmark datasets show that our model 1) outperforms previous diffusion prior based BFR methods in terms of authenticity and fidelity, and 2) outperforms the naive degradation model in modeling the real-world degradations.
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Modeling human-scene interactions (HSI) is essential for understanding and simulating everyday human behaviors. Recent approaches utilizing generative modeling have made progress in this domain; however, they are limited in controllability and flexibility for real-world applications. To address these challenges, we propose reformulating the HSI modeling problem as Scene-aware Motion In-betweening---a more tractable and practical task. We introduce SceneMI, a framework that supports several practical applications, including keyframe-guided character animation in 3D scenes and enhancing the motion quality of imperfect HSI data. SceneMI employs dual scene descriptors to comprehensively encode global and local scene context. Furthermore, our framework leverages the inherent denoising nature of diffusion models to generalize on noisy keyframes. Experimental results demonstrate SceneMI's effectiveness in scene-aware keyframe in-betweening and generalization to the real-world GIMO dataset, where motions and scenes are acquired by noisy IMU sensors and smartphones. We further showcase SceneMI's applicability in HSI reconstruction from monocular videos.
Bridging the Skeleton-Text Modality Gap: Diffusion-Powered Modality Alignment for Zero-shot Skeleton-based Action Recognition
PDF ↗In zero-shot skeleton-based action recognition (ZSAR), aligning skeleton features with the text features of action labels is essential for accurately predicting unseen actions. ZSAR faces a fundamental challenge in bridging the modality gap between the two-kind features, which severely limits generalization to unseen actions. Previous methods focus on direct alignment between skeleton and text latent spaces, but the modality gaps between these spaces hinder robust generalization learning. Motivated by the success of diffusion models in multi-modal alignment (e.g., text-to-image, text-to-video), we firstly present a diffusion-based skeleton-text alignment framework for ZSAR. Our approach, Triplet Diffusion for Skeleton-Text Matching (TDSM), focuses on cross-alignment power of diffusion models rather than their generative capability. Specifically, TDSM aligns skeleton features with text prompts by incorporating text features into the reverse diffusion process, where skeleton features are denoised under text guidance, forming a unified skeleton-text latent space for robust matching. To enhance discriminative power, we introduce a triplet diffusion (TD) loss that encourages our TDSM to correct skeleton-text matches while pushing them apart for different action classes. Our TDSM significantly outperforms very recent state-of-the-art methods with significantly large margins of 2.36%-point to 13.05%-point, demonstrating superior accuracy and scalability in zero-shot settings through effective skeleton-text matching.
Can Generative Geospatial Diffusion Models Excel as Discriminative Geospatial Foundation Models?
PDF ↗Self-supervised learning (SSL) has revolutionized representation learning in Remote Sensing (RS), advancing Geospatial Foundation Models (GFMs) to leverage vast unlabeled satellite imagery for diverse downstream tasks. Currently, GFMs primarily employ objectives like contrastive learning or masked image modeling, owing to their proven success in learning transferable representations. However, generative diffusion models, which demonstrate the potential to capture multi-grained semantics essential for RS tasks during image generation, remain underexplored for discriminative applications. This prompts the question: can generative diffusion models also excel and serve as GFMs with sufficient discriminative power? In this work, we answer this question with SatDiFuser, a framework that transforms a diffusion-based generative geospatial foundation model into a powerful pretraining tool for discriminative RS. By systematically analyzing multi-stage, noise-dependent diffusion features, we develop three fusion strategies to effectively leverage these diverse representations. Extensive experiments on remote sensing benchmarks show that SatDiFuser outperforms state-of-the-art GFMs, achieving gains of up to +5.7% mIoU in semantic segmentation and +7.9% F1-score in classification, demonstrating the capacity of diffusion-based generative foundation models to rival or exceed discriminative GFMs. The source code is available at: https://github.com/yurujaja/SatDiFuser.
In spite of the remarkable potential of Latent Diffusion Models (LDMs) in image generation, the desired properties and optimal design of the autoencoders have been underexplored. In this work, we analyze the role of autoencoders in LDMs and identify three key properties: latent smoothness, perceptual compression quality, and reconstruction quality. We demonstrate that existing autoencoders fail to simultaneously satisfy all three properties, and propose Variational Masked AutoEncoders (VMAEs), taking advantage of the hierarchical features maintained by Masked AutoEncoder. We integrate VMAEs into the LDM framework, introducing Latent Diffusion Models with Masked AutoEncoders (LDMAEs). Through comprehensive experiments, we demonstrate significantly enhanced image generation quality and computational efficiency.
QuantCache: Adaptive Importance-Guided Quantization with Hierarchical Latent and Layer Caching for Video Generation
PDF ↗Recently, Diffusion Transformers (DiTs) have emerged as a dominant architecture in video generation, surpassing U-Net-based models in terms of performance. However, the enhanced capabilities of DiTs come with significant drawbacks, including increased computational and memory costs, which hinder their deployment on resource-constrained devices. Current acceleration techniques, such as quantization and cache mechanism, offer limited speedup and are often applied in isolation, failing to fully address the complexities of DiT architectures. In this paper, we propose QuantCache, a novel training-free inference acceleration framework that jointly optimizes hierarchical latent caching, adaptive importance-guided quantization, and structural redundancy-aware pruning. QuantCache achieves an end-to-end latency speedup of 6.72x on Open-Sora with minimal loss in generation quality. Extensive evaluations across multiple video generation benchmarks demonstrate the effectiveness of our method, setting a new standard for efficient DiT inference. We will release all code and models to facilitate further research.
Textured 3D morphing creates smooth and plausible interpolation sequences between two 3D objects, focusing on transitions in both shape and texture. This is important for creative applications like visual effects in filmmaking. Previous methods rely on establishing point-to-point correspondences and determining smooth deformation trajectories, which inherently restrict them to shape-only morphing on untextured, topologically aligned datasets. This restriction leads to labor-intensive preprocessing and poor generalization. To overcome these challenges, we propose a method for 3D regenerative morphing using a 3D diffusion prior. Unlike previous methods that depend on explicit correspondences and deformations, our method eliminates the additional need for obtaining correspondence and uses the 3D diffusion prior to generate morphing. Specifically, we first introduce a 3D diffusion model and interpolate the source and target information at three levels: initial noise, model parameters, and condition features. We then explore an Attention Fusion strategy to generate smoother morphing sequences. To further improve the plausibility of semantic interpolation and the generated 3D surfaces, we propose two strategies: (a) Token Reordering, where we match approximate tokens based on semantic analysis to guide implicit correspondences in the denoising process of the diffusion model, and (b) Low-Frequency Enhancement, where we enhance low-frequency signals in the tokens to improve the quality of generated surfaces. Experimental results show that our method achieves superior smoothness and plausibility in 3D morphing across diverse cross-category object pairs, offering a novel regenerative method for 3D morphing with textured representations.
General visual representations learned from web-scale datasets for robotics have achieved great success in recent years, enabling data-efficient robot learning on manipulation tasks; yet these pre-trained representations are mostly on 2D images, neglecting the inherent 3D nature of the world. However, due to the scarcity of large-scale 3D data, it is still hard to extract a universal 3D representation from web datasets. Instead, we are seeking a general visual pre-training framework that could improve all 3D representations as an alternative. Our framework, called FVP, is a novel 4D Visual Pre-training framework for real-world robot learning. FVP frames the visual pre-training objective as a next-point-cloud-prediction problem, models the prediction model as a diffusion model, and pre-trains the model on the larger public datasets directly. Across twelve real-world manipulation tasks, FVP boosts the average success rate of 3D Diffusion Policy (DP3) for these tasks by 28%. The FVP pre-trained DP3 achieves state-of-the-art performance across imitation learning methods. Moreover, the efficacy of \ours adapts across various point cloud encoders and datasets. Finally, we apply FVP to the RDT-1B, a larger Vision-Language-Action robotic model, enhancing its performance on various robot tasks. Our project page is available at: https://4d-visual-pretraining.github.io/.
In spite of recent progress, image diffusion models still produce artifacts. A common solution is to leverage the feedback provided by quality assessment systems or human annotators to optimize the model, where images are generally rated in their entirety. In this work, we believe problem-solving starts with identification, yielding the request that the model should be aware of not just the presence of defects in an image, but their specific locations. Motivated by this, we propose DiffDoctor, a two-stage pipeline to assist image diffusion models in generating fewer artifacts. Concretely, the first stage targets developing a robust artifact detector, for which we collect a dataset of over 1M flawed synthesized images and set up an efficient human-in-the-loop annotation process, incorporating a carefully designed class-balance strategy. The learned artifact detector is then involved in the second stage to optimize the diffusion model by providing pixel-level feedback. Extensive experiments on text-to-image diffusion models demonstrate the effectiveness of our artifact detector as well as the soundness of our diagnose-then-treat design.
Driven by large-scale model iterations, the inference speed and generalization ability of 3D model generation have improved significantly. However, the quality of existing methods still falls short of enabling direct use without post-processing. Common issues include insufficient texture clarity, loss of semantic information, lack of fine-grained detail, and the generation of redundant artifacts. Moreover, current approaches focus solely on producing static structures, where individual components remain non-movable, without considering functional applications in the generation process. To address these limitations, we draw inspiration from LEGO-like modular construction and decompose complex models into semantically functional components. We propose LEGO-Maker, a novel framework that reformulates the text-to-3D task into a three-stage process: target image generation, functional semantic decomposition, and multi-task 3D generation with structured fusion. Leveraging a reorganized high-quality 3D dataset, we train a Diffusion model and a semantic segmentation model tailored for 3D generation tasks. Additionally, we design a motion-driven mechanism to introduce action sequences for functionally interactive modules after model fusion. Experimental results demonstrate that, compared to existing methods, our approach significantly enhances semantic understanding, model detail quality, and text consistency while showcasing direct applicability across various scenarios.
We present a latent diffusion model for fast feed-forward 3D scene generation. Given one or more images, our model Bolt3D directly samples a 3D scene representation in less than seven seconds on a single GPU. We achieve this by leveraging powerful and scalable existing 2D diffusion network architectures to produce consistent high-fidelity 3D scene representations. To train this model, we create a large-scale multiview-consistent dataset of 3D geometry and appearance by applying state-of-the-art dense 3D reconstruction techniques to existing multiview image datasets. Compared to prior multiview generative models that require per-scene optimization for 3D reconstruction, Bolt3D reduces the inference cost by a factor of 300 times. Project website: szymanowiczs.github.io/bolt3d
Latent diffusion models have demonstrated superior performance over traditional methods in generating highly detailed and aesthetically pleasing images, which makes them widely used for various image generation and editing tasks, including outpainting. However, most LDM-based outpainting methods impose constraints on resolution and aspect ratio, often leading to the loss of local details and blurring. One way to address these issues is progressive outpainting, where the image is extended outward incrementally. However, naive progressive outpainting suffers from two key challenges: (1) difficulty in effectively capturing global context, making it hard to maintain the original context, and (2) a tendency to generate unnatural patterns. These challenges are particularly pronounced in art, where artists pre-design the composition before painting. As a result, existing methods often introduce visual inconsistencies that distract the viewer and diminish the intended artistic emphasis. To address these limitations, we propose two types of composition planning modules that enhance progressive outpainting by leveraging global structural guidance. These modules guide a pre-trained stable diffusion model to consider the overall composition, enabling realistic and contextually appropriate artwork completion without labor-intensive user prompts. Through experiments on diverse artwork images, we show the effectiveness of our proposed method both quantitatively and qualitatively.
Efficient Input-level Backdoor Defense on Text-to-Image Synthesis via Neuron Activation Variation
PDF ↗In recent years, text-to-image (T2I) diffusion models have gained significant attention for their ability to generate high-quality images reflecting text prompts. However, their growing popularity has also led to the emergence of backdoor threats, posing substantial risks. Currently, effective defense strategies against such threats are lacking due to the diversity of backdoor targets in T2I synthesis. In this paper, we propose NaviT2I, an efficient input-level backdoor defense framework against diverse T2I backdoors. Our approach is based on the new observation that trigger tokens tend to induce significant neuron activation variation in the early stage of the diffusion generation process, a phenomenon we term Early-step Activation Variation. Leveraging this insight, NaviT2I navigates T2I models to prevent malicious inputs by analyzing Neuron activation variations caused by input tokens. Extensive experiments show that NaviT2I significantly outperforms the baselines in both effectiveness and efficiency across diverse datasets, various T2I backdoors, and different model architectures including UNet and DiT. Furthermore, we show that our method remains effective under potential adaptive attacks.
Co-Painter: Fine-Grained Controllable Image Stylization via Implicit Decoupling and Adaptive Injection
PDF ↗Controllable diffusion models have been widely applied in image stylization. However, existing methods often treat the style in the reference image as a single, indivisible entity, which makes it difficult to transfer specific stylistic attributes. To address this issue, we propose a fine-grained controllable image stylization framework, Co-Painter, to decouple multiple attributes embedded in the reference image and adaptively inject it into the diffusion model. We first build a multi-condition image stylization framework based on the text-to-image generation model. Then, to drive it, we develop a fine-grained decoupling mechanism to implicitly separate the attributes from the image. Finally, we design a gated feature injection mechanism to adaptively regulate the importance of multiple attributes. To support the above procedure, we also build a dataset with fine-grained styles. It comprises nearly 48,000 image-text pairs samples. Extensive experiments demonstrate that the proposed model achieves an optimal balance between text alignment and style similarity to reference images, both in standard and fine-grained settings.
ReFlex: Text-Guided Editing of Real Images in Rectified Flow via Mid-Step Feature Extraction and Attention Adaptation
PDF ↗Rectified Flow text-to-image models surpass diffusion models in image quality and text alignment, but adapting ReFlow for real-image editing remains challenging. We propose a new real-image editing method for ReFlow by analyzing the intermediate representations of multimodal transformer blocks and identifying three key features. To extract these features from real images with sufficient structural preservation, we leverage mid-step latent, which is inverted only up to the mid-step. We then adapt attention during injection to improve editability and enhance alignment to the target text. Our method is training-free, requires no user-provided mask, and can be applied even without a source prompt. Extensive experiments on two benchmarks with nine baselines demonstrate its superior performance over prior methods, further validated by human evaluations confirming a strong user preference for our approach.
Match-cuts are powerful cinematic tools that create seamless transitions between scenes, delivering strong visual and metaphorical connections. However, crafting impactful match-cuts is a challenging and resource-intensive process that requires deliberate artistic planning throughout the production pipeline. In this work, we introduce MatchDiffusion, a training-free method that uses text-to-video diffusion models to automatically generate match-cuts. As such, MatchDiffusion is the first method for match-cut generation. Our method leverages an inherent property of diffusion models, whereby the early denoising steps determine the broad appearance of the scene, while the latter steps add details. Motivated by this property, MatchDiffusion first performs "Joint Diffusion", by initializing generation for two prompts from a shared noise sample, and following a shared denoising path for the first denoising steps.This process results in the two videos sharing structural and motion characteristics. After Joint Diffusion, we then conduct "Disjoint Diffusion", allowing the videos' denoising paths to diverge and introduce their unique details. MatchDiffusion thus yields visually coherent videos that are amenable to match-cuts. We demonstrate the effectiveness of our method through user studies and metrics, showing its potential to democratize match-cut creation.
The Dichromatic Reflection Model (DRM), a widely used physical image formation model, has been extensively applied to specular highlight removal. However, traditional DRM solvers fail to effectively recover the missing content underneath specular highlights and are prone to producing visual artifacts. Additionally, existing deep learning-based methods do not fully exploit the full set of the DRM variables. Instead, they typically learn to translate an input image into its diffuse image (and specular residue image). As a result, their performance remains somewhat limited. To overcome these limitations, we propose a neural DRM solver for specular highlight removal. Our pipeline consists of three networks: Highlight Detection Network (HDNet), Alpha-Chrom Estimation Network (ACENet), and Refinement Network (RNet). Specifically, HDNet is first used to detect specular highlights. Meanwhile, guided by multi-level contextual contrasted features from HDNet, ACENet estimates the DRM variables. These estimated variables are then used by our reconstruction models to generate the specular-free and specular residue images. Finally, we introduce RNet, a model-free network, to further refine the results. Extensive experiments on existing datasets and our collected images demonstrate that our neural solver is superior to previous methods.
Recently, autoregressive (AR) image models have demonstrated remarkable generative capabilities, positioning themselves as a compelling alternative to diffusion models. However, their sequential nature leads to long inference times, limiting their practical scalability. In this work, we introduce Grouped Speculative Decoding (GSD), a novel, training-free acceleration method for AR image models. While recent studies have explored Speculative Decoding (SD) as a means to speed up AR image generation, existing approaches either provide only modest acceleration or require additional training. Our in-depth analysis reveals a fundamental difference between language and image tokens: image tokens exhibit inherent redundancy and diversity, meaning multiple tokens can convey valid semantics. However, traditional SD methods are designed to accept only a single most-likley token, which fails to leverage this difference, leading to excessive false-negative rejections. To address this, we propose a new SD strategy that evaluates clusters of visually valid tokens rather than relying on a single target token. Additionally, we observe that static clustering based on embedding distance is ineffective, which motivates our dynamic GSD approach. Extensive experiments show that GSD accelerates AR image models by an average of 3.7x while preserving image quality--all without requiring any additional training.
PromptDresser: Improving the Quality and Controllability of Virtual Try-On via Generative Textual Prompt and Prompt-aware Mask
PDF ↗Recent virtual try-on approaches have advanced by finetuning pre-trained text-to-image diffusion models to leverage their powerful generative ability; however, the use of text prompts in virtual try-on remains underexplored. This paper tackles a text-editable virtual try-on task that modifies the clothing based on the provided clothing image while editing the wearing style (e.g., tucking style, fit) according to the text descriptions. In the text-editable virtual try-on, three key aspects exist: (i) designing rich text descriptions for paired person-clothing data to train the model, (ii) addressing the conflicts where textual information of the existing person's clothing interferes the generation of the new clothing, and (iii) adaptively adjust the inpainting mask aligned with the text descriptions, ensuring proper editing areas while preserving the original person's appearance irrelevant to the new clothing. To address these aspects, we propose PromptDresser, a text-editable virtual try-on model that leverages large multimodal model (LMM) assistance to enable high-quality and versatile manipulation based on generative text prompts. Our approach utilizes LMMs via in-context learning to generate detailed text descriptions for person and clothing images independently, including pose details and editing attributes using minimal human cost. Moreover, to ensure the editing areas, we adjust the inpainting mask depending on the text prompts adaptively. Our approach enhances text editability while effectively conveying clothing details that are difficult to capture through images alone, leading to improved image quality. Experiments show that PromptDressersignificantly outperforms baselines, demonstrating superior text-driven control and versatile clothing manipulation.
MotionDiff: Training-free Zero-shot Interactive Motion Editing via Flow-assisted Multi-view Diffusion
PDF ↗Generative models have made remarkable advancements and are capable of producing high-quality content. However, performing controllable editing with generative models remains challenging, due to their inherent uncertainty in outputs. This challenge is particularly pronounced in motion editing, which involves the processing of spatial information. While some physics-based generative methods have attempted to implement motion editing, they typically operate on single-view images with simple motions, such as translation and dragging. These methods struggle to handle complex rotation and stretching motions and ensure multi-view consistency, often necessitating resource-intensive retraining. To address these challenges, we propose a training-free zero-shot diffusion method that leverages optical flow for complex multi-view motion editing. Specifically, given a static scene, users can interactively select objects of interest to add motion priors. The proposed Point Kinematic Model then estimates corresponding multi-view optical flows during the Multi-view Flow Estimation Stage. Subsequently, these optical flows are utilized to generate multi-view motion results through decoupled motion representation in the Multi-view Motion Diffusion Stage. Extensive experiments demonstrate that our method outperforms other physics-based generative motion editing methods in achieving high-quality multi-view consistent motion results. Notably, our method does not require retraining, enabling users to conveniently adapt it for various down-stream tasks.