This work presents AnyDoor a diffusion-based image generator with the power to teleport target objects to new scenes at user-specified locations with desired shapes. Instead of tuning parameters for each object our model is trained only once and effortlessly generalizes to diverse object-scene combinations at the inference stage. Such a challenging zero-shot setting requires an adequate characterization of a certain object. To this end we complement the commonly used identity feature with detail features which are carefully designed to maintain appearance details yet allow versatile local variations(e.g. lighting orientation posture etc.) supporting the object in favorably blending with different surroundings. We further propose to borrow knowledge from video datasets where we can observe various forms (i.e. along the time axis) of a single object leading to stronger model generalizability and robustness. Extensive experiments demonstrate the superiority of our approach over existing alternatives as well as its great potential in real-world applications such as virtual try-on shape editing and object swapping.
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Humans can infer 3D structure from 2D images of an object based on past experience and improve their 3D understanding as they see more images. Inspired by this behavior we introduce SAP3D a system for 3D reconstruction and novel view synthesis from an arbitrary number of unposed images. Given a few unposed images of an object we adapt a pre-trained view-conditioned diffusion model together with the camera poses of the images via test-time fine-tuning. The adapted diffusion model and the obtained camera poses are then utilized as instance-specific priors for 3D reconstruction and novel view synthesis. We show that as the number of input images increases the performance of our approach improves bridging the gap between optimization-based prior-less 3D reconstruction methods and single-image-to-3D diffusion-based methods. We demonstrate our system on real images as well as standard synthetic benchmarks. Our ablation studies confirm that this adaption behavior is key for more accurate 3D understanding.
Personalization has emerged as a prominent aspect within the field of generative AI enabling the synthesis of individuals in diverse contexts and styles while retaining high-fidelity to their identities. However the process of personalization presents inherent challenges in terms of time and memory requirements. Fine-tuning each personalized model needs considerable GPU time investment and storing a personalized model per subject can be demanding in terms of storage capacity. To overcome these challenges we propose HyperDreamBooth - a hypernetwork capable of efficiently generating a small set of personalized weights from a single image of a person. By composing these weights into the diffusion model coupled with fast finetuning HyperDreamBooth can generate a person's face in various contexts and styles with high subject details while also preserving the model's crucial knowledge of diverse styles and semantic modifications. Our method achieves personalization on faces in roughly 20 seconds 25x faster than DreamBooth and 125x faster than Textual Inversion using as few as one reference image with the same quality and style diversity as DreamBooth. Also our method yields a model that is 10000x smaller than a normal DreamBooth model.
In this paper we present a method to reconstruct the world and multiple dynamic humans in 3D from a monocular video input. As a key idea we represent both the world and multiple humans via the recently emerging 3D Gaussian Splatting (3D-GS) representation enabling to conveniently and efficiently compose and render them together. In particular we address the scenarios with severely limited and sparse observations in 3D human reconstruction a common challenge encountered in the real world. To tackle this challenge we introduce a novel approach to optimize the 3D-GS representation in a canonical space by fusing the sparse cues in the common space where we leverage a pre-trained 2D diffusion model to synthesize unseen views while keeping the consistency with the observed 2D appearances. We demonstrate our method can reconstruct high-quality animatable 3D humans in various challenging examples in the presence of occlusion image crops few-shot and extremely sparse observations. After reconstruction our method is capable of not only rendering the scene in any novel views at arbitrary time instances but also editing the 3D scene by removing individual humans or applying different motions for each human. Through various experiments we demonstrate the quality and efficiency of our methods over alternative existing approaches.
Text-guided image editing finds applications in various creative and practical fields. While recent studies in image generation have advanced the field they often struggle with the dual challenges of coherent image transformation and context preservation. In response our work introduces prompt augmentation a method amplifying a single input prompt into several target prompts strengthening textual context and enabling localised image editing. Specifically we use the augmented prompts to delineate the intended manipulation area. We propose a Contrastive Loss tailored to driving effective image editing by displacing edited areas and drawing preserved regions closer. Acknowledging the continuous nature of image manipulations we further refine our approach by incorporating the similarity concept creating a Soft Contrastive Loss. The new losses are incorporated to the diffusion model demonstrating improved or competitive image editing results on public datasets and generated images over state-of-the-art approaches.
We present Condition-Aware Neural Network (CAN) a new method for adding control to image generative models. In parallel to prior conditional control methods CAN controls the image generation process by dynamically manipulating the weight of the neural network. This is achieved by introducing a condition-aware weight generation module that generates conditional weight for convolution/linear layers based on the input condition. We test CAN on class-conditional image generation on ImageNet and text-to-image generation on COCO. CAN consistently delivers significant improvements for diffusion transformer models including DiT and UViT. In particular CAN combined with EfficientViT (CaT) achieves 2.78 FID on ImageNet 512x512 surpassing DiT-XL/2 while requiring 52x fewer MACs per sampling step.
We present personalized residuals and localized attention-guided sampling for efficient concept-driven generation using text-to-image diffusion models. Our method first represents concepts by freezing the weights of a pretrained text-conditioned diffusion model and learning low-rank residuals for a small subset of the model's layers. The residual-based approach then directly enables application of our proposed sampling technique which applies the learned residuals only in areas where the concept is localized via cross-attention and applies the original diffusion weights in all other regions. Localized sampling therefore combines the learned identity of the concept with the existing generative prior of the underlying diffusion model. We show that personalized residuals effectively capture the identity of a concept in 3 minutes on a single GPU without the use of regularization images and with fewer parameters than previous models and localized sampling allows using the original model as strong prior for large parts of the image.
LeftRefill: Filling Right Canvas based on Left Reference through Generalized Text-to-Image Diffusion Model
PDF ↗This paper introduces LeftRefill an innovative approach to efficiently harness large Text-to-Image (T2I) diffusion models for reference-guided image synthesis. As the name implies LeftRefill horizontally stitches reference and target views together as a whole input. The reference image occupies the left side while the target canvas is positioned on the right. Then LeftRefill paints the right-side target canvas based on the left-side reference and specific task instructions. Such a task formulation shares some similarities with contextual inpainting akin to the actions of a human painter. This novel formulation efficiently learns both structural and textured correspondence between reference and target without other image encoders or adapters. We inject task and view information through cross-attention modules in T2I models and further exhibit multi-view reference ability via the re-arranged self-attention modules. These enable LeftRefill to perform consistent generation as a generalized model without requiring test-time fine-tuning or model modifications. Thus LeftRefill can be seen as a simple yet unified framework to address reference-guided synthesis. As an exemplar we leverage LeftRefill to address two different challenges: reference-guided inpainting and novel view synthesis based on the pre-trained StableDiffusion. Codes and models are released at https://github.com/ewrfcas/LeftRefill.
While AI-generated content has garnered significant attention achieving photo-realistic video synthesis remains a formidable challenge. Despite the promising advances in diffusion models for video generation quality the complex model architecture and substantial computational demands for both training and inference create a significant gap between these models and real-world applications. This paper presents SNED a superposition network architecture search method for efficient video diffusion model. Our method employs a supernet training paradigm that targets various model cost and resolution options using a weight-sharing method. Moreover we propose the supernet training sampling warm-up for fast training optimization. To showcase the flexibility of our method we conduct experiments involving both pixel-space and latent-space video diffusion models. The results demonstrate that our framework consistently produces comparable results across different model options with high efficiency. According to the experiment for the pixel-space video diffusion model we can achieve consistent video generation results simultaneously across 64 x 64 to 256 x 256 resolutions with a large range of model sizes from 640M to 1.6B number of parameters for pixel-space video diffusion models.
KTPFormer: Kinematics and Trajectory Prior Knowledge-Enhanced Transformer for 3D Human Pose Estimation
PDF ↗This paper presents a novel Kinematics and Trajectory Prior Knowledge-Enhanced Transformer (KTPFormer) which overcomes the weakness in existing transformer-based methods for 3D human pose estimation that the derivation of Q K V vectors in their self-attention mechanisms are all based on simple linear mapping. We propose two prior attention modules namely Kinematics Prior Attention (KPA) and Trajectory Prior Attention (TPA) to take advantage of the known anatomical structure of the human body and motion trajectory information to facilitate effective learning of global dependencies and features in the multi-head self-attention. KPA models kinematic relationships in the human body by constructing a topology of kinematics while TPA builds a trajectory topology to learn the information of joint motion trajectory across frames. Yielding Q K V vectors with prior knowledge the two modules enable KTPFormer to model both spatial and temporal correlations simultaneously. Extensive experiments on three benchmarks (Human3.6M MPI-INF-3DHP and HumanEva) show that KTPFormer achieves superior performance in comparison to state-of-the-art methods. More importantly our KPA and TPA modules have lightweight plug-and-play designs and can be integrated into various transformer-based networks (i.e. diffusion-based) to improve the performance with only a very small increase in the computational overhead. The code is available at: https://github.com/JihuaPeng/KTPFormer.
Classifier-Free Guidance (CFG) has been widely used in text-to-image diffusion models where the CFG scale is introduced to control the strength of text guidance on the whole image space. However we argue that a global CFG scale results in spatial inconsistency on varying semantic strengths and suboptimal image quality. To address this problem we present a novel approach Semantic-aware Classifier-Free Guidance (S-CFG) to customize the guidance degrees for different semantic units in text-to-image diffusion models. Specifically we first design a training-free semantic segmentation method to partition the latent image into relatively independent semantic regions at each denoising step. In particular the cross-attention map in the denoising U-net backbone is renormalized for assigning each patch to the corresponding token while the self-attention map is used to complete the semantic regions. Then to balance the amplification of diverse semantic units we adaptively adjust the CFG scales across different semantic regions to rescale the text guidance degrees into a uniform level. Finally extensive experiments demonstrate the superiority of S-CFG over the original CFG strategy on various text-to-image diffusion models without requiring any extra training cost. our codes are available at https://github.com/SmilesDZgk/S-CFG.
We present a simple yet effective technique to estimate lighting in a single input image. Current techniques rely heavily on HDR panorama datasets to train neural networks to regress an input with limited field-of-view to a full environment map. However these approaches often struggle with real-world uncontrolled settings due to the limited diversity and size of their datasets. To address this problem we leverage diffusion models trained on billions of standard images to render a chrome ball into the input image. Despite its simplicity this task remains challenging: the diffusion models often insert incorrect or inconsistent objects and cannot readily generate chrome balls in HDR format. Our research uncovers a surprising relationship between the appearance of chrome balls and the initial diffusion noise map which we utilize to consistently generate high-quality chrome balls. We further fine-tune an LDR diffusion model (Stable Diffusion XL) with LoRA enabling it to perform exposure bracketing for HDR light estimation. Our method produces convincing light estimates across diverse settings and demonstrates superior generalization to in-the-wild scenarios.
Text-to-image (T2I) diffusion models notably the unCLIP models (e.g. DALL-E-2) achieve state-of-the-art (SOTA) performance on various compositional T2I benchmarks at the cost of significant computational resources. The unCLIP stack comprises T2I prior and diffusion image decoder. The T2I prior model alone adds a billion parameters compared to the Latent Diffusion Models which increases the computational and high-quality data requirements. We introduce ECLIPSE a novel contrastive learning method that is both parameter and data-efficient. ECLIPSE leverages pre-trained vision-language models (e.g. CLIP) to distill the knowledge into the prior model. We demonstrate that the ECLIPSE trained prior with only 3.3% of the parameters and trained on a mere 2.8% of the data surpasses the baseline T2I priors with an average of 71.6% preference score under resource-limited setting. It also attains performance on par with SOTA big models achieving an average of 63.36% preference score in terms of the ability to follow the text compositions. Extensive experiments on two unCLIP diffusion image decoders Karlo and Kandinsky affirm that ECLIPSE priors consistently deliver high performance while significantly reducing resource dependency. Project page: https://eclipse-t2i.vercel.app/
Text2QR: Harmonizing Aesthetic Customization and Scanning Robustness for Text-Guided QR Code Generation
PDF ↗In the digital era QR codes serve as a linchpin connecting virtual and physical realms. Their pervasive integration across various applications highlights the demand for aesthetically pleasing codes without compromised scannability. However prevailing methods grapple with the intrinsic challenge of balancing customization and scannability. Notably stable-diffusion models have ushered in an epoch of high-quality customizable content generation. This paper introduces Text2QR a pioneering approach leveraging these advancements to address a fundamental challenge: concurrently achieving user-defined aesthetics and scanning robustness. To ensure stable generation of aesthetic QR codes we introduce the QR Aesthetic Blueprint (QAB) module generating a blueprint image exerting control over the entire generation process. Subsequently the Scannability Enhancing Latent Refinement (SELR) process refines the output iteratively in the latent space enhancing scanning robustness. This approach harnesses the potent generation capabilities of stable-diffusion models navigating the trade-off between image aesthetics and QR code scannability. Our experiments demonstrate the seamless fusion of visual appeal with the practical utility of aesthetic QR codes markedly outperforming prior methods. Codes are available at https://github.com/mulns/Text2QR
MicroDiffusion: Implicit Representation-Guided Diffusion for 3D Reconstruction from Limited 2D Microscopy Projections
PDF ↗Volumetric optical microscopy using non-diffracting beams enables rapid imaging of 3D volumes by projecting them axially to 2D images but lacks crucial depth information. Addressing this we introduce MicroDiffusion a pioneering tool facilitating high-quality depth-resolved 3D volume reconstruction from limited 2D projections. While existing Implicit Neural Representation (INR) models often yield incomplete outputs and Denoising Diffusion Probabilistic Models (DDPM) excel at capturing details our method integrates INR's structural coherence with DDPM's fine-detail enhancement capabilities. We pretrain an INR model to transform 2D axially-projected images into a preliminary 3D volume. This pretrained INR acts as a global prior guiding DDPM's generative process through a linear interpolation between INR outputs and noise inputs. This strategy enriches the diffusion process with structured 3D information enhancing detail and reducing noise in localized 2D images.By conditioning the diffusion model on the closest 2D projection MicroDiffusion substantially enhances fidelity in resulting 3D reconstructions surpassing INR and standard DDPM outputs with unparalleled image quality and structural fidelity. Our code and dataset are available athttps://github.com/UCSC-VLAA/MicroDiffusion.
Physically-Based Rendering (PBR) is key to modeling the interaction between light and materials and finds extensive applications across computer graphics domains. However acquiring PBR materials is costly and requires special apparatus. In this paper we propose a method to extract PBR materials from a single real-world image. We do so in two steps: first we map regions of the image to material concept tokens using a diffusion model allowing the sampling of texture images resembling each material in the scene. Second we leverage a separate network to decompose the generated textures into spatially varying BRDFs (SVBRDFs) offering us readily usable materials for rendering applications. Our approach relies on existing synthetic material libraries with SVBRDF ground truth. It exploits a diffusion-generated RGB texture dataset to allow generalization to new samples using unsupervised domain adaptation (UDA). Our contributions are thoroughly evaluated on synthetic and real-world datasets. We further demonstrate the applicability of our method for editing 3D scenes with materials estimated from real photographs. Along with video we share code and models as open-source on the project page: https://github.com/astra-vision/MaterialPalette
HyperSDFusion: Bridging Hierarchical Structures in Language and Geometry for Enhanced 3D Text2Shape Generation
PDF ↗3D shape generation from text is a fundamental task in 3D representation learning. The text-shape pairs exhibit a hierarchical structure where a general text like "chair" covers all 3D shapes of the chair while more detailed prompts refer to more specific shapes. Furthermore both text and 3D shapes are inherently hierarchical structures. However existing Text2Shape methods such as SDFusion do not exploit that. In this work we propose HyperSDFusion a dual-branch diffusion model that generates 3D shapes from a given text. Since hyperbolic space is suitable for handling hierarchical data we propose to learn the hierarchical representations of text and 3D shapes in hyperbolic space. First we introduce a hyperbolic text-image encoder to learn the sequential and multi-modal hierarchical features of text in hyperbolic space. In addition we design a hyperbolic text-graph convolution module to learn the hierarchical features of text in hyperbolic space. In order to fully utilize these text features we introduce a dual-branch structure to embed text features in 3D feature space. At last to endow the generated 3D shapes with a hierarchical structure we devise a hyperbolic hierarchical loss. Our method is the first to explore the hyperbolic hierarchical representation for text-to-shape generation. Experimental results on the existing text-to-shape paired dataset Text2Shape achieved state-of-the-art results. We release our implementation under HyperSDFusion.github.io.
We present Intrinsic Image Diffusion a generative model for appearance decomposition of indoor scenes. Given a single input view we sample multiple possible material explanations represented as albedo roughness and metallic maps. Appearance decomposition poses a considerable challenge in computer vision due to the inherent ambiguity between lighting and material properties and the lack of real datasets. To address this issue we advocate for a probabilistic formulation where instead of attempting to directly predict the true material properties we employ a conditional generative model to sample from the solution space. Furthermore we show that utilizing the strong learned prior of recent diffusion models trained on large-scale real-world images can be adapted to material estimation and highly improves the generalization to real images. Our method produces significantly sharper more consistent and more detailed materials outperforming state-of-the-art methods by 1.5dB on PSNR and by 45% better FID score on albedo prediction. We demonstrate the effectiveness of our approach through experiments on both synthetic and real-world datasets.
We present MM-AU a novel dataset for Multi-Modal Accident video Understanding. MM-AU contains 11727 in-the-wild ego-view accident videos each with temporally aligned text descriptions. We annotate over 2.23 million object boxes and 58650 pairs of video-based accident reasons covering 58 accident categories. MM-AU supports various accident understanding tasks particularly multimodal video diffusion to understand accident cause-effect chains for safe driving. With MM-AU we present an Abductive accident Video understanding framework for Safe Driving perception (AdVersa-SD). AdVersa-SD performs video diffusion via an Object-Centric Video Diffusion (OAVD) method which is driven by an abductive CLIP model. This model involves a contrastive interaction loss to learn the pair co-occurrence of normal near-accident accident frames with the corresponding text descriptions such as accident reasons prevention advice and accident categories. OAVD enforces the object region learning while fixing the content of the original frame background in video generation to find the dominant objects for certain accidents. Extensive experiments verify the abductive ability of AdVersa-SD and the superiority of OAVD against the state-of-the-art diffusion models. Additionally we provide careful benchmark evaluations for object detection and accident reason answering since AdVersa-SD relies on precise object and accident reason information.
SIFU: Side-view Conditioned Implicit Function for Real-world Usable Clothed Human Reconstruction
PDF ↗Creating high-quality 3D models of clothed humans from single images for real-world applications is crucial. Despite recent advancements accurately reconstructing humans in complex poses or with loose clothing from in-the-wild images along with predicting textures for unseen areas remains a significant challenge. A key limitation of previous methods is their insufficient prior guidance in transitioning from 2D to 3D and in texture prediction. In response we introduce SIFU (Side-view Conditioned Implicit Function for Real-world Usable Clothed Human Reconstruction) a novel approach combining a Side-view Decoupling Transformer with a 3D Consistent Texture Refinement pipeline. SIFU employs a cross-attention mechanism within the transformer using SMPL-X normals as queries to effectively decouple side-view features in the process of mapping 2D features to 3D. This method not only improves the precision of the 3D models but also their robustness especially when SMPL-X estimates are not perfect. Our texture refinement process leverages text-to-image diffusion-based prior to generate realistic and consistent textures for invisible views. Through extensive experiments SIFU surpasses SOTA methods in both geometry and texture reconstruction showcasing enhanced robustness in complex scenarios and achieving an unprecedented Chamfer and P2S measurement. Our approach extends to practical applications such as 3D printing and scene building demonstrating its broad utility in real-world scenarios.